falText to imageComing soon

fal-ai/fooocus

Name: fal-ai/fooocus
Brand: fal

fal-ai/fooocus is a text to image tool from fal that runs in the background and returns images. Coming soon: listed with its price, not yet open to runs.

This tool is listed but not yet available to run.

Price pending

Pricing

Price pending.

Limits

Runs: In the background
Returns: image
Longest queue wait: 3600s
Run deadline: 1800s
Largest inline input: 1024 KB
Largest input file: 500 MB
Largest output: 4768 MB
Outputs per run: 16

Input

Field	Kind	Default	Description
`prompt`	text		The prompt to use for generating the image. Be as descriptive as possible for best results.
`negative_prompt`	text		The negative prompt to use. Use it to address details that you don't want in the image. This could be colors, objects, scenery and even the small details (e.g. moustache, blurry, low resolution).
`styles`	list	["Fooocus Sharp","Fooocus V2","Fooocus Enhance"]	The style to use.
`performance`	one of Speed, Quality, Extreme Speed, Lightning	Extreme Speed	You can choose Speed or Quality
`guidance_scale`	number 1 to 30	4	The CFG (Classifier Free Guidance) scale is a measure of how close you want the model to stick to your prompt when looking for a related image to show you.
`sharpness`	number 0 to 30	2	The sharpness of the generated image. Use it to control how sharp the generated image should be. Higher value means image and texture are sharper.
`aspect_ratio`	text	1024x1024	The size of the generated image. You can choose between some presets or custom height and width that must be multiples of 8.
`num_images`	number (whole) 1 to 4	1	Number of images to generate in one request
`loras`	JSON	[{"path":"https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/resolve/main/sd_xl_offset_example-lora_1.0.safetensors","scale":0.1}]	The LoRAs to use for the image generation. You can use up to 5 LoRAs and they will be merged together to generate the final image.
`refiner_model`	one of None, realisticVisionV60B1_v51VAE.safetensors	None	Refiner (SDXL or SD 1.5)
`refiner_switch`	number 0 to 1	0.8	Use 0.4 for SD1.5 realistic models; 0.667 for SD1.5 anime models 0.8 for XL-refiners; or any value for switching two SDXL models.
`output_format`	one of png, jpeg, webp	jpeg	The format of the generated image.
`seed`	number (whole)	—	The same seed and the same prompt given to the same version of Stable Diffusion will output the same image every time.
`control_image_url`	text	—	The image to use as a reference for the generated image.
`control_type`	one of ImagePrompt, PyraCanny, CPDS, FaceSwap	PyraCanny	The type of image control
`control_image_weight`	number 0 to 2	1	The strength of the control image. Use it to control how much the generated image should look like the control image.
`control_image_stop_at`	number 0 to 1	1	The stop at value of the control image. Use it to control how much the generated image should look like the control image.
`inpaint_image_url`	text	—	The image to use as a reference for inpainting.
`mask_image_url`	text	—	The image to use as a mask for the generated image.
`mixing_image_prompt_and_inpaint`	true or false	false

Same request in code

curl https://api.spiral.ai/v1/runs \
  -H "Authorization: Bearer $SPIRAL_API_KEY" \
  -H "Content-Type: application/json" \
  -H "Idempotency-Key: $(uuidgen)" \
  -d '{
    "provider": "fal",
    "model": "fal-ai/fooocus",
    "input": {}
  }'